Stable Diffusion

Stable Diffusion is an open latent text-to-image diffusion architecture offering versatile fine-tuning (LoRA, ControlNet), local inference, and precise prompt typography adherence.

When to Use

• Open-Source Generative Image Synthesis: Generating photorealistic or stylized artwork using SDXL and Flux with Hugging Face diffusers.
• Precise Spatial Conditioning with ControlNet: Guiding image composition with depth maps, canny edges, and human poses.
• Inpainting, Outpainting & Image-to-Image: Modifying specific masked regions while preserving surrounding image context.
• Custom Style Adaptation with LoRA: Applying fine-tuned community style and character adapters.

Quick Start

import torch
from diffusers import StableDiffusionXLPipeline